G1Joshi

stable-diffusion

Stable Diffusion image generation models. Use for image AI.

G1Joshi 9 2 Updated 3mo ago
GitHub

Install

npx skillscat add g1joshi/agent-skills/stable-diffusion

Install via the SkillsCat registry.

SKILL.md

Stable Diffusion

Stable Diffusion (by Stability AI) is the open standard for image generation. SD 3.5 (2025) improves prompt adherence and typography.

When to Use

  • Control: You need exact composition control (ControlNet).
  • Local Generation: Run on your own GPU. No censorship/cost.
  • Fine-Tuning: training LoRAs on your own face/product.

Core Concepts

Diffusers Library

The Hugging Face library to run SD pipelines in Python.

LoRA (Low-Rank Adaptation)

Small adapter files (100MB) that add a style or character to the base model.

ComfyUI

The node-based GUI for building complex SD workflows (see separate skill).

Best Practices (2025)

Do:

  • Use SD 3.5 Large: For best text rendering.
  • Use Flux: The community has largely moved to Flux.1 (Black Forest Labs) alongside SD.
  • Use ControlNet: To force the image to follow a specific pose or edge map.

Don't:

  • Don't use SD 1.5: Unless you need specific legacy LoRAs. SDXL/SD3/Flux are superior.

References